Popular GitHub Repositories for Huggingface Diffusers

Create ? videos with Stable Diffusion by exploring the latent space and morphing between text prompts

Dreambooth implementation based on Stable Diffusion with minimal code.

? Unofficial huggingface/diffusers-based implementation of the paper "Training-Free Structured Diffusion Guidance for Compositional Text-to-Image Synthesis".

A simple web application that lets you replace any part of an image with an image generated based on your description.

Official repository of "TryOffDiff: Virtual-Try-Off via High-Fidelity Garment Reconstruction using Diffusion Models".

A fine-tuned model based on Stable Diffusion to create images in the style of Midjourney

Quantized stable-diffusion cutting down memory 75%, testing in streamlit, deploying in container

? HuggingFace Diffusers Flax TPU and PyTorch GPU for Colab

Dreambooth for colab

Morpheus is an open-source project that offers a creative and innovative platform for generating stunning artworks using image editing and stable diffusion models.

Diffusers API in OCaml

Easily create your own AI avatar images!

This is a repo providing same stable diffusion experiments, regarding textual inversion task and captioning task

diffusion model for unconditional image generation of Bored Apes

Configuration files for building E621-Rising v3 SDXL model and dataset

Experimental Stable Diffusion XL Webui

Generating images with diffusion models on a mobile device, with an intranet GPU box as backend

扩散模型训练利器：Diffusers 与 PyTorch-Lightning 强强联手，打造简洁灵活的训练模板。

Generate image by prompts using Stable Diffusion XL-turbo, in Pytorch & Gradio

Here I have used Stable Diffusion with the diffusers from hugging face, took one image as input and then added elements into it by using diffusion algorithm, and iterated this process three/four times with different type of elements, adding by text inputs.